Stable Diffusion XL 1.0

Stable Diffusion XL 1.0 is Stability AI's foundational text-to-image model. With a 3.5B parameter ensemble (base + refiner), it generates 1024×1024 images and was the model that kicked off the open image generation revolution. Still one of the most downloaded models at 1.74M monthly.

Architecture

SDXL uses a two-stage pipeline:

1. Base model (3.5B) - generates latents from text prompt
2. Refiner model - enhances detail and quality

Usage

from diffusers import DiffusionPipeline
import torch
 
pipe = DiffusionPipeline.from_pretrained(
    "stabilityai/stable-diffusion-xl-base-1.0",
    torch_dtype=torch.float16
)
pipe.to("cuda")
 
image = pipe(
    prompt="a photo of a cat on a park bench",
    height=1024, width=1024
).images[0]

When to use

• Legacy workflows - huge ecosystem of LoRAs and extensions
• Creative art - flexible and well-understood
• Open research - fully open weights (CreativeML Open RAIL-M)